generate_image_from_image
Transform an existing image into a new one based on a prompt, with adjustable strength to preserve structure. Apply style transfer, refine details, or maintain composition.
Instructions
Generate an image using img2img (image-to-image) transformation.
Takes an existing image and transforms it based on the prompt while preserving structure according to the strength parameter.
Use cases:
Style transfer (apply anime/comic/realistic style to photo)
Image refinement (improve details, fix artifacts)
Pose/composition preservation (keep layout, change style)
Character consistency (transform existing character art)
Input Schema
| Name | Required | Description | Default |
|---|---|---|---|
| seed | No | Random seed for reproducibility (random if not set) | |
| loras | No | List of LoRA filenames to apply (in order) | |
| steps | No | Inference steps (default 28, use 4-8 with LCM LoRA) | |
| prompt | Yes | The positive prompt describing desired output | |
| strength | No | Denoising strength (0.0-1.0). Controls how much to change. - 0.0-0.3: Minor refinements, preserve most details - 0.3-0.5: Moderate changes, good for style transfer - 0.5-0.7: Significant changes, keeps composition - 0.7-1.0: Major transformation, only basic structure preserved | |
| checkpoint | No | Checkpoint filename or 'default' for HuggingFace model | |
| image_path | Yes | Absolute path to source image to transform | |
| lora_scales | No | Scale/strength per LoRA (0.0-2.0, defaults to 1.0) | |
| render_type | No | Optional render type specification ('gpu' or 'cpu'). If specified and doesn't match detected device, renders are aborted to prevent slow processing. | |
| guidance_scale | No | CFG scale (default 5.0) | |
| negative_prompt | No | Optional; defaults to standard negative prompt |
Output Schema
| Name | Required | Description | Default |
|---|---|---|---|
No arguments | |||