generate_image
Generate images from text prompts using a diffusion model with customizable checkpoints, LoRA styles, resolution, and seed. Outputs include metadata for reproducible workflows.
Instructions
Generate an image with Animagine XL 4.0.
Uses the Diffusers pipeline with the lpw_stable_diffusion_xl custom pipeline. Images are saved to outputs/YYYY-MM-DD/ with accompanying metadata JSON.
Supports custom checkpoints and LoRA mixing for style control.
Recommended workflow:
list_models → see available checkpoints and LoRAs
validate_prompt → check for issues
optimize_prompt → improve structure
generate_image → create the image
Input Schema
| Name | Required | Description | Default |
|---|---|---|---|
| seed | No | Random seed for reproducibility (random if not set) | |
| loras | No | List of LoRA filenames to apply (in order). Examples: ["custom_lora.safetensors"] (user's local LoRAs) | |
| steps | No | Inference steps (default 28, use 4-8 with LCM LoRA) | |
| width | No | Image width (default 832, portrait) | |
| height | No | Image height (default 1216, portrait) | |
| prompt | Yes | The positive prompt (pre-validated recommended) | |
| checkpoint | No | Checkpoint filename or 'default' for HuggingFace model. Examples: "custom_checkpoint.safetensors" (user's local checkpoint) | |
| lora_scales | No | Scale/strength per LoRA (0.0-2.0, defaults to 1.0 for each). Example: [0.8, 0.5] for two LoRAs | |
| render_type | No | Optional render type specification ('gpu' or 'cpu'). If specified and doesn't match detected device, renders are aborted to prevent slow processing. | |
| guidance_scale | No | CFG scale (default 5.0, use 1.5 with LCM LoRA) | |
| negative_prompt | No | Optional; defaults to standard negative prompt |
Output Schema
| Name | Required | Description | Default |
|---|---|---|---|
No arguments | |||