generate
Create images from text prompts on a local Stable Diffusion Forge. For image-to-image edits, supply an init image and denoising strength to control how closely the output follows the original.
Instructions
Generate an image from an already-written prompt.
The prompt is sent verbatim: include any <lora:name:weight> yourself. With
use_profile_defaults on, missing sampling parameters are filled from what the
loaded model actually used before, so leave them unset unless you mean to
override. That includes shift (Forge's distilled_cfg_scale) and the
dimensions: leaving them at 0 takes the architecture's own values instead of
a generic default. Returns file paths when the output folder is readable.
Pass init_image (a local file path) to run img2img instead, where
denoising_strength controls how far the result may drift from it: around
0.3 keeps the composition, 0.75 reinterprets it freely. Edit-style and video
models expect values close to 1.0.
Input Schema
| Name | Required | Description | Default |
|---|---|---|---|
| seed | No | ||
| shift | No | ||
| steps | No | ||
| width | No | ||
| height | No | ||
| prompt | Yes | ||
| cfg_scale | No | ||
| scheduler | No | ||
| batch_size | No | ||
| init_image | No | ||
| sampler_name | No | ||
| negative_prompt | No | ||
| denoising_strength | No | ||
| use_profile_defaults | No |